Linear Algebra II: Eigenvectors and Diagonalisability

This post continues the discussion of the Oxford first-year course Linear Algebra II. We’ve moved on from determinants, and are now considering eigenvalues and eigenvectors of matrices and linear maps.

A good question to ask is: what’s the point of knowing about eigenvectors? I can think of a quick answer and a longer answer. The quick answer is that whenever we have a mapping of any kind, it is natural to ask about its fixed points. And since we are thinking about vector spaces and linear maps, if we can’t find any fixed points, we might nonetheless be able to find the best thing, some vectors whose direction is fixed by the map. In general, knowing about fixed points of a mapping might tell us other more qualitative properties, including the behaviour seen when you apply the map iteratively a large number of times. (Indeed a recent post discusses this exact problem for positive matrices in a context relevant to a chapter of my thesis…)

A more specific answer concerns bases. Recall that a linear map is defined independently of any basis: it’s just a map from the vector space to itself. We can express the linear map via a matrix with respect to some basis, but how to choose the basis? We could always choose the canonical basis in \mathbb{R}^n, since it’s easy to do vector and matrix calculations when most of the entries of all the vectors are zero. We also have a good visual idea (at least in up to three dimensions) of what a matrix might mean with respect to that basis. If we needed to divide the three-dimensional world around us into small volumes, we’d tend to describe it with small cubes rather than small arbitrary parallelopipeds.

But once we know something about the linear map, we might want to choose a basis of vectors on which the behaviour of the map is particularly easy to describe. And eigenvectors fulfil precisely this role. If we are able to choose a basis of eigenvectors, describing the map’s action, either abstractly, or via a (diagonal) matrix, is very straightforward. If we are given a matrix to begin with, we know how to do a change of basis, and changing to the basis of eigenvectors is precisely what’s going when we write A=P^{-1}DP, where D is a diagonal matrix. We construct P by taking its columns to be these eigenvectors. In particular, for a given vector x, y=Px is the vector giving the coefficients of x in the basis of eigenvectors.

So the case where we have a basis of eigenvectors is particularly useful, and in this case, we say the matrix or the map is diagonalisable. Remember how we find eigenvalues. If there exists a non-zero vector x satisfying Ax=\lambda x, then x is in the kernel of A-\lambda I. As we discussed last time, introducing the determinant gives a much more manageable way to verify which values of \lambda result in A-\lambda I having a non-trivial kernel. In particular, if non-zero x is in the kernel, we have \mathrm{det}(A-\lambda I)=0, and this leads to a polynomial of degree n (the dimensional of the vector space / size of the matrix) for \lambda, called the characteristic polynomial \chi_A(z), which has the eigenvalues as its roots.

If we agree to work over the complex field, then this is good, because it means we always have eigenvalues, and so it becomes sensible to talk about exactly how many eigenvalues and eigenvectors we have. Observe that if we restrict to real vector spaces, this might not be the case. In the plane, the rotation by \pi/2 for example has no fixed vectors.

Multiplicities of eigenvalues

We call the algebraic multiplicity \alpha(\lambda) of an eigenvalue \lambda to be the exponent of the factor (z-\lambda) in the factorisation of the characteristic polynomial. To define the geometric multiplicity, observe that all the eigenvectors with eigenvalue \lambda form a subspace, and so it is meaningful to talk about the dimension of this subspace (‘eigenspace’), which is the geometric multiplicity \gamma(\lambda). There are two facts that one needs to remember. The slightly less obvious one is that \gamma(\lambda)\le \alpha(\lambda) for all \lambda. One can see this by, for example, working in a basis that extends a basis of the \lambda-eigenspace. Observe at this stage that the sum of the algebraic multiplicities has to be n by definition, while the sum of geometric multiplicities is at most n. And this makes sense, because the space spanned by all the eigenvectors is a subspace, and so has dimension at most n.

The more obvious, but more frequently forgotten result is that

\alpha(\lambda)\ge 1 \quad \iff \quad \gamma(\lambda)\ge 1,

which is simply a consequence of the property discussed a few paragraphs previously concerning the kernel of A-\lambda I.

In particular, we might make the heuristic observation that ‘most’ polynomials of degree n have n distinct roots. This is certainly true for quadratics: there is only one value that the discriminant can take such that we see a repeated root. Alternatively, imagine shifting the quadratic up and down (in a complex way if necessary); again there is only one moment at which it might have a repeated root. This observation can be generalised easily to higher degree polynomials in a number of ways.

So if we lift this observation across to matrices, we see that most matrices have n distinct eigenvalues, and thus have n linearly independent eigenvectors which form a basis, hence the matrix is diagonalisable. I think it’s really worth reflecting on this, since much of a first exploration into linear algebra ends up treating exactly the case where the matrix is not diagonalisable.

The principal example of a non-diagonalisable matrix is \begin{pmatrix}2&1\\0&2\end{pmatrix}, where the 2s can be replaced by any value, and the 1 can be replaced by an non-zero value. There’s plenty to learn about to what extent versions of this matrix of higher size represent all non-diagonalisable matrices, but such an exposition of Jordan normal form comes next year for the students taking this course.

It probably is worth saying now though, that this example gives a good sanity check for whether a method is actually using diagonalisability correctly. For example, it is easily seen that elementary row operations to not preserve diagonalisability by starting from \begin{pmatrix}2&0\\0&2\end{pmatrix} and ending up at our counter-example. One could also argue from this that the set of non-diagonalisable matrices are dense within the set of matrices with a repeated eigenvalue. That is, having a repeated eigenvalue but full eigenspace is doubly-infinitely-unlikely.

Cayley-Hamilton theorem

Anyway, among other results, we also saw the Cayley-Hamilton theorem, which states that a matrix A satisfies its own characteristic equation. That is \chi_A(A)=0, where the zero on the right-hand side is the zero matrix. It’s tempting to substitute A into the expression \mathrm{det}(A-\lambda I), but of course this is not valid. Indeed imagine a typical eigenvalue determinant matrix with terms like (7-\lambda) on the diagonal; it doesn’t make sense to substitute a matrix for \lambda as one of the entries of the overall matrix!

Fortunately, we can argue convincingly in the case where A is a diagonalisable matrix. Remember that \chi_A(A) is a matrix. Now looki at the action of \chi_A(A) on any eigenvector v, corresponding to eigenvalue \lambda. Applying some power of A to v gives v multiplied by the same power of \lambda, and so we end up with

\chi_A(A)v = \chi_A(\lambda)v = 0.

This only worked when v was an eigenvector, but fortunately there is a basis of eigenvectors if A is diagonalisable, and so \chi_A(A)v=0 for all v, hence \chi_A(A)=0.

But \chi_A(A) is just a matrix-valued function of A. If you think about it, \chi_A is a monic polynomial, all of whose non-leading coefficients are multinomials of degree at most n-1 in the entries of A. Furthermore, these multinomials have (non-negative) integer coefficients. Therefore the entries of \chi_A(A) are multinomials of degree at most 2n-1 in the entries of A, and again have (non-negative) integer coefficients.

Even without the integrality of the coefficients, this says that, under any reasonable definition of continuity of matrices (which could be induced from any topology on \mathbb{R}^{n\times n}) the function \chi_A(A) should be continuous as a function of A. But we’ve shown \chi_A(A)=0 for all diagonalisable A, and also argued that most complex-valued matrices are diagonalisable. Turning this into a formal statement about denseness means that we’ve shown the Cayley-Hamilton theorem for non-diagonalisable matrices also. It feels that because the coefficients are non-negative integers, we might also have shown the result for other fields too, but I have minimal knowledge or recollection at the moment of the things one has to check for this sort of result.

It’s worth ending with the brief comment that Cayley-Hamilton is useful, among other reasons because it enables us to write the inverse of A as a polynomial of degree at most n-1 in terms of A. In many settings this is a lot easier to work with in terms of calculations than an argument with minors.


Linear Algebra II: Determinants 2

In the previous post, we introduced determinants of matrices (and by extension linear maps) via its multilinearity properties, and as the change-of-volume factor. We also discussed how to calculate them, via row operations, or Laplace expansion, or directly via a sum of products of entries over permutations.

The question of why this is ever a useful quantity to consider remains, and this post tries to answer it. We’ll start by seeing one example where this is a very natural quantity to consider, and then the main abstract setting, where the determinant is zero, and consider a particularly nice example of this.

Jacobeans as a determinant

We consider integration by substitution. Firstly, in one variable: when it comes to Riemann integration of a function g(x) with respect to x, we view dx as the width of a small column which approximates the function near x. Now, if we reparameterise, that is if we write x=f(y) for some well-behaved (in particular differentiable) function f, then the width of the column is dx= dy.(dx/dy)=f'(y) dy. This may be negative, if y is decreasing while x is increasing, but for now let’s not worry about this overly, for example by assuming the function g is non-negative. Thus if we want to integrate with y as the variable, we multiply the integrand by this factor |f'(y)|.

What about in higher dimensions? We have exactly the same situation, only instead of two-dimensional columns, we have (n+1)-dimensional columns. We then multiply the n-dimensional volume of the base by the height, again given by g(\mathbf{x}). If we have a similar transformation of the base variable \mathbf{x}=f(\mathbf{y}), we differentiate to get

\mathrm{d}x_i = \sum_{j=1}^n\frac{\mathrm{d}f_i}{\mathrm{d}y_j} \mathrm{d}y_j.

In other words

\mathrm{d}\mathbf{x}= J \mathrm{d}\mathbf{y},

where J is the Jacobean matrix of partial deriatives. In particular, we know how to relate the volume [0,\mathrm{d}x_1]\times\ldots\times [0,\mathrm{d}x_n] to the volume [0,\mathrm{d}y_1]\times \ldots\times [0,\mathrm{d}y_n]. It’s simply the determinant of the Jacobean J. So if we want to integrate with respect to \mathbf{y}, it only remains to pre-multiply the integrand by |\mathrm{det}J| and proceed otherwise as in the one-dimensional case.

Det A = 0

A first linear algebra course might well motivate the introducing matrices as a notational shortcut for solving families of linear equations, Ax=b. The main idea is that generally we can solve this equation uniquely. Almost all of the theory developed in such a first linear algebra course deals with the case when this fails to hold. In particular, there are many ways to characterise this case, and we list some of them now:

  • Ax=b has no solutions for some b;
  • A is not invertible;
  • A has non-trivial kernel, that is, with dimension at least one;
  • A does not have full rank, that is, the image has dimension less than n;
  • The columns (or indeed the rows) are linearly dependent;
  • The matrix can be row-reduced to a matrix with a row of zeroes.

It is useful that these are equivalent, as in abstract problems one can choose whichever interpretation from this list is most relevant. However, all of these are quite hard to check. Exhibiting a non-trivial kernel element is hard – one either has to do manual row-reduction, or the equivalent in the context of linear equations. But we can add the characterisation

  • det A = 0;

to the list. And this is genuinely much easier to check for specific examples, either abstract or numerical.

Let’s quickly convince ourselves of a couple of these equivalences. Determinant is invariant under row-reductions, and by multilinearity it is certainly the case that det A = 0 if A has a row of zeroes. We also said that A is the change-of-volume factor. Note that A is a map from the domain to its image, so if A has less than full rank, then any set in the image has zero volume.

The Vandermonde matrix

This is a good example of this theory in practice. Consider the Vandermonde matrix where each row is a geometric progression:

V=\begin{pmatrix}1&\alpha_1&\ldots&\alpha_1^{n-1}\\1&\alpha_2&\ldots&\alpha_2^{n-1}\\ \vdots&\vdots&\ddots&\vdots\\ 1&\alpha_n&\ldots&\alpha_n^{n-1}\end{pmatrix}.

Now suppose we attempt to solve

V\begin{pmatrix}a_0\\a_1\\ \vdots\\ a_{n-1}\end{pmatrix}=\begin{pmatrix}b_1\\b_2\\ \vdots \\ b_n\end{pmatrix}.

There’s a natural interpretation to this, that’s especially clear with this suggestive notation. Each row corresponds to a polynomial, where the coefficients are given by the (a_0,a_1,\ldots,a_{n-1}), and the argument is given by \alpha_i.

So if we try to solve for (a_0,a_1,\ldots,a_{n-1}), given (\alpha_1,\ldots,\alpha_n) and (b_1,\ldots,b_n), we are asking whether we can find a polynomial P with degree at most n-1 such that P(\alpha_i)=b_i for $i=1. Lagrange interpolation gives an argument where we just directly write down the relevant polynomial, but we can also deploy our linear algebraic arguments too.

The equivalence of all these statements means that to verify existence and uniqueness of such a polynomial, we only need to check that the Vandermonde matrix has non-zero determinant. And in fact there are a variety of methods to show that

\mathrm{det}V=\prod_{1\le i< j}\le n(\alpha_j-\alpha_i).

For the polynomial question to be meaningful, we would certainly demand that the (\alpha_i) are distinct, and so this determinant is non-zero, and we’ve shown that n points determine a degree (n-1) polynomial uniquely.

If we multiply on the left instead, suppose that we are considering a discrete probability distribution X that takes n known values (\alpha_1,\ldots,\alpha_n) with unknown probabilities (p_1,\ldots,p_n). Then we have

(p_1,\ldots,p_n) V = (1,\mathbb{E}X, \mathbb{E}[X^2],\ldots, \mathbb{E}[X^{n-1}]).

So, again by inverting the Vandermonde matrix (which is know is possible since its determinant is non-zero…) we can recover the distribution from the first (n-1) moments of the distribution.

A similar argument applies to show that the Discrete Fourier Transform is invertible, and in this case (where the \alpha_is are roots of unity), the expression for the Vandermonde determinant is particularly tractable.

Linear Algebra II: Determinants 1

This term, I’m giving tutorials on a course that’s new to me, the apparently notorious ‘Linear Algebra II’ for first year undergraduates. I can appreciate how it might have ended up with this reputation, but as always, every challenge is also an opportunity, So I’m going to (try to) write a short series of blog posts about what we’ve discussed in the tutorials.

The first problem sheet-and-a-half concerned determinants of matrices. There are three things worth addressing here:

  1. What are abstract definitions, and which is most useful in each setting?
  2. How to actually calculate them?
  3. What’s the overall point?

The answers are obviously not completely unrelated, but we’ll probably defer the third question to a second post.

The determinant is a map from the set of matrices \mathcal{M}_n to the base field (hereafter assumed to be \mathbb{R},\mathbb{C}). The Oxford course defines it through its properties:

  • Multilinear in the columns of the matrix.
  • Equal to zero if two columns are equal.
  • Equal to one if the matrix is the identity.

One then checks that there is a unique such map, and so from now on it’s reasonable to call it the determinant of the matrix. It will follow from pretty much any consequence that we can replace ‘columns’ with ‘rows’ throughout and get the same map.

Other definitions

We have a closed form expression for the determinant given via permutations of n

\mathrm{det}(A)=\sum_{\sigma\in \Sigma_n} \mathrm{sign}(\sigma) a_{1\sigma(1)}\ldots a_{n\sigma(n)}.

We’ll come back to a discussion of when this particular definition is useful. It can be derived by carefully transforming the identity matrix into A, using the operations which are mentioned in the original definition of the determinant, in particular, keeping track of the number of transpositions of columns.

It’s clear from any definition that the determinant is a polynomial of degree n in the entries of the matrix, but this definition will be useful if you want to make some more precise comment on the nature of this polynomial. For example, if entries of the matrix are polynomials in x of various degree (think of the eigenvalue equation for example) this allows you to control (or at least bound) the overall degree of the determinant as a polynomial in x.

The determinant is also the volume of the n-dimensional parallelopiped formed by the column vectors of the matrix. This is easy to check in two dimensions, for the matrix \begin{pmatrix}a&b\\c&d\end{pmatrix}:

20160225_172206To calculate the area of the central parallelogram, we have to subtract the area of two small rectangles and four small triangles from the outer rectangle, obtaining

(a+b)(c+d)-2bc - \frac12(ac+ac+bd+bd)=ad-bc,

as we expect. This calculation is harder to execute in higher dimensions, and certainly harder to visualise.

Maybe, though, we don’t have to, so long as we can reassure ourselves that this volume satisfies the implicit definition of the determinant map at the start. Multilinearity in the columns is not that hard to see. If we multiply the jth column by some constant, we are stretching the parallelopiped by the same constant factor in one direction, and so the volume grows appropriately. The additivity property can similarly be thought of as joining together two parallelopipeds at their common face (which is common since the other column vectors have to be constant in this construction). If two column vectors are equal, then clearly this volume actually has dimension at most n-1, and thus volume zero, so the final two conditions are genuinely easy to check.

The challenge here is that there is a direction involved. Determinants can be negative, but in our classical viewpoint, areas generally are not. In 2D, we can think of this as saying that the area is positive if the vector (b,d) lies anti-clockwise from (a,c) in the parallelogram, while is it negative otherwise. Again, this is harder to visualise in higher dimensions, but it is at least plausible that one could develop a similar decomposition. Ultimately, we are happy with the notion of directed lengths (ie vectors on the real line), and these are easy to add up without having to separate into cases, and the case holds for areas and higher-dimensional volumes.

Evaluating determinants

If we actually want to compute the determinant of a given matrix, the sum over permutations is intractable since it doesn’t have any natural splits into stages. The implicit definitions and this area consideration are clearly useless for all but the most special of examples.

The Laplace expansion is the usual algorithm to calculate the determinant of an n x n matrix. You pick a row (or a column), and evaluate the determinants of the (n-1) x (n-1) minor matrices given by deleting this row (or column), and each column (or row) in turn. This leaves us with n determinants of smaller matrices, which we pre-multiply by the entries in the original deleted row (or column), and add up in an alternating way (*). This is highly computationally intensive for large matrices, but for 3×3 and 4×4 can be done by hand with probability of an error bounded away from 1.

There is the flexibility to choose the reference row or column. Since the entries of these affect the sum through small products, it is highly convenient to choose a row or column with a lot of zeros. In particular, if there’s a row or column with exactly one non-zero entry, this is an ideal candidate.

The sum over permutations also works well when a lot of the entries are zero, because then a lot of the permutations give a summand which is zero. Upper-triangular matrices are a good example: only one permutation (the identity permutation) avoids all the zero elements underneath the diagonal.

One can also observe from the multilinearity property of the determinant map that there are lots of operations we can apply to the matrix which leave the determinant fixed. These are often called elementary row operations, though obviously we can apply them to the columns as well. To summarise, if we interchange two rows, the sign of the determinant is reversed. And if we add some multiple of one row to any other row, the determinant stays the same.

When matrices are not square, it’s quite important to be specific about exactly what form you can reduce a general matrix to via such row operations, but in this context, it’s not hugely important. Reduced echelon form (without the condition that leading coefficients will be one) is achievable, but this is a special case of an upper triangular matrix, for which the determinant is given by the product of the diagonal entries, ie is easy.

Whether this is substantially easier than Laplace expansion depends on the matrix itself and taste, both to do manually, and to code.

(*) I’m not a fan personally of this alternating definition. It seems to me much more natural to define the minor as

M^{i,j}=(a_{i+k,j_\ell})_{1\le k,\ell\le n-1},

with indices taken modulo n. Then you don’t have any \pm 1s in the Laplace expansion.

Using determinants in abstract problems

So the determinant gives directly the area of the image (under A) of the unit hypercube. By linearity (of A), it is easy to see that it also gives the scale factor of the area change (under A) of any hyper-cuboid, parallel to the conventional axes, anywhere in the space. Then, eg by approximating any sensible n-dimensional shape (*) as a union of such hyper-cuboids, we can show that in fact the area of any sensible shape increases by a factor (det A) under application of A.

This is a good thing to remember, because it is an excellent heuristic for seeing why the determinant of a linear map is basis-independent. It also gives a much easier proof of the key result

\mathrm{det}(AB)=\mathrm{det}(A) \mathrm{det}(B),

than that given by fixing B and viewing det(AB) as a map from matrices to the field, just like the original definition of determinant.

Some of the theory in the course is proved using elementary row operations. But these invite complicated notation, so are best used only in simple arguments, or when things are fairly explicit to begin with. Given an abstract problem about determinants of matrices, it is often tempting to induct on the size of the matrix in some way. I think it’s worth saying that even though the Laplace expansion is explicitly set up in this way, the notation involved is also likely to be annoying here, while permutations are easy to describe inductively: eg let \sigma(1)=k, then view the remainder of the permutation as a bijection \{2,3,\ldots,n\}\rightarrow [n]\backslash \{k\}.

Shortly, we’ll have a second post answering the final question: what’s the point of working with determinants? We’ve already seen half an answer, in that they describe the change-of-volume factor of a matrix (or linear map), but this can be substantially developed.

Characterisations of Eigenvalues

I’ve been working for much of the past few months on a version of the frozen percolation random graph process with types. The connectivity between types is controlled by a (finite) non-negative square matrix, and so I’ve been engaging with linear algebra theory to an extent I haven’t really experienced since the second or third year of undergraduate maths.

We are interested in whether the graphs in question are subcritical, critical or supercritical. As in the case of multitype branching processes, this is controlled by the principal eigenvalue of a related non-negative matrix. So I’ve been looking up lots of methods for controlling eigenvalues, and some have proved useful, and some have not, but I thought it would be worthwhile to present some of them here.

Bounds and characterisations of spectral radius

Throughout, I will be talking about finite, square matricies. Eigenvalues may be defined as roots of the characteristic polynomial, and so by the fundamental theorem of algebra, there is always at least one complex eigenvalue. There is always at least one eigenvector associated to any eigenvalue. However, the dimension of the eigenspace is not always the same as the multiplicity of the eigenvalue as a root of the characteristic polynomial. The latter is called algebraic multiplicity, while the former is geometric multiplicity.

For now though, this distinction will be unimportant. The spectral radius of a matrix A is defined as

\rho(A)=\max \{|\lambda|\, : \, \lambda \text{ and eigenvalue of }A\}.

We can bound the spectral radius in terms of the norm of the matrix. Remember that a matrix norm has to satisfy all the usual properties of a norm, as well as a submultiplicative property |||AB|||\le |||A|||\cdot |||B|||. This is good, as otherwise we would be free to replace any norm by an arbitrary multiple of itself, and so no useful bounds could ever emerge. Note that the submultiplicativity implies that |||I_n||\ge 1.

Now, let \lambda,x be some eigenvalue and associated (right-)eigenvector respectively of matrix A. Let X be the square matrix given by taking all the columns to be x. Now Ax=\lambda x implies AX=\lambda X, and so

|\lambda| \cdot|||X||| = |||\lambda X||| = |||AX||| \le |||A|||\cdot |||X|||,

and thus we conclude our most basic bound \lambda \le |||A|||.

When A is diagonalisable, life is particularly easy, but in general we can write A as a conjugate of its Jordan normal form. Then, by looking at each diagonal block of the Jordan normal form separately, we can show that

\lim_{k\rightarrow 0}A^k = 0\quad \iff \quad \rho(A)<1.

Then, applying this, with additional care, to the matrices A / (\rho(A)\pm \epsilon), we derive Gelfand’s Formula, that \rho(A) = \lim_{k\rightarrow \infty} ||A^k||^{1/k}. Again, this applies for any matrix norm.

Real symmetric matrices

When the matrix is real and symmetric, it is not too hard to show that all the eigenvalues are real, and furthermore that all the geometric multiplicities are equal to the algebraic multiplicities. That is, the matrix is diagonalisable, and there is an (orthogonal) basis of eigenvectors. Once we assume we are working with respect to this eigenbasis, it is easy to see how the Rayleigh quotient characterisation of the largest (and smallest) eigenvalue works. Let’s say the eigenvalues are \lambda_1\ge \lambda_2\ge\ldots \ge \lambda_n, then for any ||x||_2=1, we have \lambda_1\ge x^T A x\ge \lambda_n, and equality is attained when x is the respective eigenvector, normalised appropriately.

This is an especially useful characterisation of the largest eigenvalue, as for example we can see fairly easily that this means \lambda_1 is a convex function of the (real, symmetric) matrix.

We can generalise this Rayleigh quotient idea if we take k orthonormal vectors in R^k, arrange them in an nxk matrix P, so that P^T P = I_k. Now we consider the matrix P^TAP. [Note that if k=1, we are exactly considering x^TA x as before.] Then Poincare’s Separation Theorem say that the eigenvalues \mu_1\ge \mu_2\ge\ldots \mu \mu_k of P^TAP (which is also real, symmetric) are bounded by the original eigenvalues:

\lambda_{n-k+i} \ge \mu_i\ge \lambda_i.

Since the trace is preserved under conjugation, and the trace is the sum of eigenvalues, we can apply this result with P’s columns taken to be the any k canonical basis vectors of \mathbb{R}^k. Without loss of generality, we may assume the basis has been chosen so that the diagonal elements of A satisfy a_{11}\ge a_{22}\ge\ldots\ge a_{nn}, and so now we have that the sequence (a_{11},a_{22},\ldots,a_{nn}) is majorised by (\lambda_1,\lambda_2,\ldots,\lambda_n) and majorises (\lambda_n,\lambda_{n-1},\ldots,\lambda_1). The first of these relations can be used via the setup of Karamata’s inequality to conclude that for any convex function f, we have

\sum_{i=1}^n f(\lambda_i)\ge \sum_{i=1}^n f(a_{ii}).

Gershgorin Circles

In fact, we can relate the eigenvalues to the diagonal entries of the matrix in a more general setting. We are motivated by the thought that if the off-diagonal entries are all very small, then the set of eigenvalues should be approximately given by the set of diagonal entries.

For a square complex matrix, let \lambda,x be an eigenvalue, eigenvector pair. For any index i, we have

\lambda - a_{ii}= \frac{\sum_j a_{ij}x_j}{x_i} - a_{ii} = \frac{\sum_{j\ne i}a_{ij}x_j}{x_i}.

Now consider the i such that x_i=\max |x_j|, and take absolute values and apply the triangle inequality,

|\lambda - a_{ii}| \le \sum_{j\ne i} \left| \frac{a_{ij}x_j}{x_i} \right| \le \sum_{j\ne i}|a_{ij}|.

Let’s define R_i=\sum_{j\ne i}|a_{ij}| to be the sum of the non-diagonal entries of the ith row. Then the Gershgorin circle theorem says that every eigenvalue lies within at least one of the discs B(a_{ii},R_i), in the complex plane. So our motivation still makes sense. If the off-diagonal entries are small, this is a strong restriction, and if they are not typically smaller than the diagonal entries, then we perhaps do not learn very much. Obviously, we could apply the same argument to the columns too.

When the diagonal entries are distinct, and the off-diagonal entries are small, the Gershgorin discs are distinct, and we would expect each to contain exactly one eigenvalue, corresponding to the appropriate diagonal entry. In fact, we can say something stronger. In general, the union of the discs is a subset of the complex plane with some connected components. Then, if a component is the union of exactly r discs, then it contains exactly r of the eigenvalues.

To see this, consider multiplying all the off-diagonal entries by z\in[0,1] and observe what happens as z varies from 0 to 1. When z=0, the matrix is diagonal, and each eigenvalue is in the Gershgorin disc (which is a single complex number). As z varies continuously, the characteristic polynomial varies continuously, and also its roots, that is the set of eigenvalues. So since each of the r eigenvalues are initially within the union of the r original, large Gershgorin discs, they must remain within this union as z varies, since they cannot ‘jump’ to another component.

It’s hard to know how time will allow, but provisionally in the next post I will talk about how to control the evolution of eigenvectors as a function of the matrix, and in particular what can go wrong.


For the middle section, I used the progression from Chapter 4 of Matrix Differential Calculus with Applications in Statistics and Econometrics (Magnus and Neudecker).

Invariant Distributions of Markov Chains

My lecture course in Linyi was all about Markov chains, and we spent much of the final two sessions discussing the properties of invariant distributions. I was not surprised, however, that none of the class chose this topic as the subject for a presentation to give after the end of the teaching week. One of the main problems is that so many rather similar properties are introduced roughly simultaneously. As we did in the class, I thought it was worth making some sort of executive summary, as a mixture of revision and interest.

Definition: \pi is an invariant measure if \pi P=\pi. If in addition \sum_{i\in I}\pi_i=1, then we say it is an invariant distribution. Of course, if I is finite, then any invariant measure can be normalised to give an invariant distribution.

The key initial questions are about existence and uniqueness. First, if there are multiple communicating classes, then an invariant measure (resp. distribution) is a linear (resp. affine) combination of the invariant measures / distributions on each (closed) class. So we restrict attention to irreducible Markov chains.

In the finite case, P is a stochastic matrix so has a column eigenvector with eigenvalue 1, namely the vector with all entries equal to 1. Thus, by reference to general theory in linear algebra, P has a row eigenvector \pi with eigenvalue 1. To paraphrase a remark made by one of my students, what is not clear is that this should be a measure. Demonstrating that this is true is rather non-trivial I think, normally done by reference to the rather more general Perron-Frobenius theorem, though on the flight home I came up with a short argument using Lagrangian duality. For now, we accept existence in the finite case, and note that we typically show existence by showing that the vector of expected time spent in each state between successive visits to a fixed reference state satisfies the properties of an invariant measure.

This is a good moment to note that recurrence is not a necessary condition for the existence of an invariant measure. For example, the random walk on \mathbb{Z}^3 is transient, but the uniform measure is invariant. However, it is not a sufficient condition for the existence of an invariant distribution either. (Of course, an irreducible finite chain is always recurrent, and always has an invariant distribution, so now we are considering only the infinite state space case.) The random walk on \mathbb{Z}^2 is recurrent, but the invariant measure is not normalisable.

The property we in fact need is positive recurrence. This says that the expected return time to each point is finite. Again, this is a class property. This is a common requirement in probabilistic arguments: almost surely finite is often not strong enough to show results if the expectation is infinite (see for example the various requirements for the optional stopping theorem). If this holds, then \pi_i=\frac{1}{\mathbb{E}T_i}, where T_i is the the return time starting from some i\in I.

The final question is ‘Why are we interested?’ One of the best answers is to look at convergence properties. A simple suggestion is this: if we start in equilibrium, then X_0,X_1,X_2,\ldots are all equal in distribution. Note that the dependence structure remains complicated, and much much more interesting than the individual distributions. Next, we observe that a calculation of n-step transition probabilities for a finite chain will typically involve a linear combination of nth powers of eigenvalues. One of the eigenvalues is 1, and the others lie strictly between -1 and 1. We observe in examples that the constant coefficient in p_{ij}^{(n)} is generally a function of j alone, and so p_{ij}^{(n)}\rightarrow\lambda_j, some distribution on I. By considering P^{n+1}=P\cdot P^n, it is easy to see that if this converges, (\lambda_j)_{j\in I} is an invariant distribution. The classic examples which do not work are

P=\begin{pmatrix}0&1\\1&0\end{pmatrix} and P=\begin{pmatrix}0&1&0\\ 0&0&1\\1&0&0\end{pmatrix},

as then the distribution of X_n is a function of the remainder of n modulo 3 alone. With a little thought, we can give a precise classification of such chains which force you to be in particular proper subsets of the state space at regular times n. Chains without this property are called aperiodic, and we can show that distributions for such chains converge to the equilibrium distribution as n\rightarrow\infty.